The Best 2238 Image Generation Tools in 2025

Stable Diffusion is a latent text-to-image diffusion model capable of generating realistic images from any text input.

A text-to-image generation model based on stable diffusion with image inpainting capabilities

SDXL 1.0 is a diffusion-based text-to-image generation model that employs an expert-integrated latent diffusion process, supporting high-resolution image generation

Stable Diffusion is a latent text-to-image diffusion model capable of generating realistic images from any text input.

The SD-XL 1.0 Refiner Model is an image generation model developed by Stability AI, designed to enhance the quality of images generated by the SDXL base model, with particular expertise in the final denoising step.

A diffusion-based text-to-image generation model that supports image generation and modification through text prompts

A latent text-to-image diffusion model based on Stable Diffusion XL, capable of image inpainting via masks

Diffusion-based text-to-image model capable of generating high-quality images from text prompts

An open-source text-to-image model capable of generating aesthetic images at 1024x1024 resolution and various aspect ratios, leading in aesthetic quality within the open-source domain.

SD-Turbo is a high-speed text-to-image model capable of generating realistic images from text prompts with just a single network inference. Released as a research prototype, it aims to explore compact distilled text-to-image models.

Animagine XL 3.0 is the latest open-source anime text-to-image model developed based on Stable Diffusion XL, focusing on high-quality anime image generation with improved hand anatomy and conceptual understanding.

Animagine XL 3.1 is an upgraded version of the Animagine XL V3 series. It has been optimized on the basis of the previous work and can generate higher-quality anime images.

IP-Adapter-FaceID is a face recognition-based text-to-image generation model capable of generating various styles of images based on facial ID embeddings.

Emi is an AI dedicated to AI art image generation, developed using Stable Diffusion XL 1.0 technology, without training on unauthorized images

A text-to-image generation model based on latent diffusion model with image inpainting capability

A diffusion model-based text-to-image generation system supporting high-resolution image generation

InstructPix2Pix is a model capable of editing images based on textual instructions, learning to translate natural language commands into image editing operations.

A text-to-image generation model fine-tuned from Stable Diffusion v2-base, utilizing latent diffusion technology and OpenCLIP text encoder

An anime-style text-to-image generation model fine-tuned based on Stable Diffusion XL

ControlNet is a neural network architecture that controls Stable Diffusion image generation through Canny edge detection conditions

The Latent Consistency Model is a text-to-image generation model distilled from a fine-tuned version of Stable-Diffusion v1-5's Dreamshaper v7, capable of producing high-quality images in extremely short inference times.

An enhanced version of the Stable Diffusion XL model, specializing in high-precision photographic style output

All-round image generation and editing control network, supporting 12 control conditions and 5 advanced editing functions

A photorealistic LoRA checkpoint trained for the FLUX.1-dev model to enhance the realism of generated images

GGUF quantized version of the FLUX.1-dev model, focused on text-to-image and image generation tasks

A multi-style text-to-image model fine-tuned on Stable Diffusion v1.5, excelling in artistic and anime styles

This is a pixel art image generation model based on Stable Diffusion, crafted with 184 training images and 8,000 training steps, specifically optimized for facial detail representation.

Hyper-SD is an advanced diffusion model acceleration technology that supports fast inference for various base models (such as FLUX.1-dev, SD3, SDXL, and SD1.5).

This is a LoRA adapter designed for Stable Diffusion v1-5, reducing inference steps to just 2-8 steps, significantly improving generation speed.

Photon is a text-to-image generation model based on stable diffusion technology, focusing on high-quality image generation.

A Stable Diffusion model fine-tuned on the animation style of Studio Ghibli, capable of generating images with Ghibli's aesthetic characteristics

Pixart-α is an efficient text-to-image generation model based on Transformer architecture, capable of generating high-quality 1024-pixel images at extremely low training costs

A text-to-image generation model based on stable diffusion technology, capable of producing high-quality, detail-rich images from textual descriptions.

A Stable Diffusion XL model specialized in illustration generation, developed based on Kohaku XL Beta 5 and trained by the OnomaAI research team

InstantID is an advanced, tuning-free method that achieves identity-preserving generation with just a single image, supporting multiple downstream tasks.

SDXL-Lightning is an ultra-fast text-to-image model that can generate high-quality 1024px images in just a few steps

Animagine XL 4.0 is an ultimate anime-themed text-to-image model fine-tuned based on Stable Diffusion XL 1.0. It is trained with 8.4 million anime-style images, significantly improving output stability and image quality.

PhotoMaker V2 is a text-to-image generation model that enables customized portrait or artwork creation without any training.

DreamShaper is a text-to-image generation tool based on the Stable Diffusion model, specializing in artistic creation, particularly excelling in generating anime-style images.

A LoRA model focused on generating amateur photography style images, suitable for diffusion model frameworks

INT4 quantized version of FLUX.1-Fill-dev, capable of filling regions in existing images based on text descriptions, providing approximately 4x memory savings and 2-3x speed improvements.

A text-to-image model trained on 124,000 Midjourney v4 images using the Stable Diffusion v1.5 framework